Dreamshaper 3.3

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from any text input.

A text-to-image generation model based on stable diffusion with image inpainting capabilities

SDXL 1.0 is a diffusion-based text-to-image generation model that employs an expert-integrated latent diffusion process, supporting high-resolution image generation

Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from any text input.

The SD-XL 1.0 Refiner Model is an image generation model developed by Stability AI, designed to enhance the quality of images generated by the SDXL base model, with particular expertise in the final denoising step.

A diffusion-based text-to-image generation model that supports image generation and modification through text prompts

A latent text-to-image diffusion model based on Stable Diffusion XL, capable of image inpainting via masks

Diffusion-based text-to-image model capable of generating high-quality images from text prompts

An open-source text-to-image model capable of generating aesthetic images at 1024x1024 resolution and various aspect ratios, leading in aesthetic quality within the open-source domain.

SD-Turbo is a high-speed text-to-image model capable of generating realistic images from text prompts with just a single network inference. Released as a research prototype, it aims to explore compact distilled text-to-image models.