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Stable Diffusion V1 5

Developed by stable-diffusion-v1-5
Stable Diffusion is a latent text-to-image diffusion model capable of generating realistic images from any text input.
Downloads 3.7M
Release Time : 8/30/2024

Model Overview

A diffusion-based text-to-image generation model using CLIP ViT-L/14 as the text encoder, primarily used for generating and modifying images based on text prompts.

Model Features

High-Quality Image Generation
Capable of generating realistic 512x512 resolution images from text prompts
Classifier-Free Guidance Sampling
Improves sampling by dropping 10% of text conditions
Multi-Framework Support
Supports various usage methods including Diffusers library, ComfyUI, Automatic1111

Model Capabilities

Text-to-Image Generation
Image Modification
Artistic Creation

Use Cases

Artistic Creation
Concept Art Generation
Generate concept artworks based on textual descriptions
Useful for rapid prototyping in fields like gaming and film
Education & Research
Generative Model Research
Used to explore and understand the limitations and biases of generative models
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