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Stable Diffusion Xl 1.0 Inpainting 0.1

Developed by diffusers
A latent text-to-image diffusion model based on Stable Diffusion XL, capable of image inpainting via masks
Downloads 673.14k
Release Time : 9/1/2023

Model Overview

This model can generate realistic images from text inputs and supports local image inpainting through masks. Initialized with stable-diffusion-xl-base-1.0 weights and trained at 1024x1024 resolution.

Model Features

High-resolution image inpainting
Supports inpainting tasks for high-resolution images up to 1024x1024
Text-guided inpainting
Guides image inpainting content based on text prompts
Synthetic mask training
Generates synthetic masks during training, with 25% cases involving fully masked images

Model Capabilities

Text-to-image generation
Image inpainting
Local image editing

Use Cases

Creative Design
Artwork restoration
Restoring old or damaged artworks
Repairs damaged areas while preserving the original style
Creative image editing
Modifying local image content based on text prompts
For example, replacing objects on a bench with a tiger
Education & Research
Generative model research
Investigating the limitations and possibilities of image inpainting techniques
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