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Stable Diffusion V1.5

Developed by stablediffusiontutorials
A latent diffusion model for text-to-image generation, supporting 512x512 resolution image generation
Downloads 1,291
Release Time : 9/8/2024

Model Overview

Stable Diffusion is a latent diffusion model capable of generating realistic images from text prompts, utilizing CLIP ViT-L/14 text encoder, suitable for artistic creation, design assistance, and other scenarios

Model Features

High-Resolution Generation
Supports 512x512 resolution image generation with improved output quality over earlier versions
Classifier-Free Guidance
Optimizes sampling process with 10% text condition dropout technique
Commercial-Friendly License
Allows commercial use and service deployment under responsible AI terms
Multi-Framework Support
Supports both Diffusers library and original GitHub repository usage

Model Capabilities

Text-to-image generation
Image style transfer
Artistic creation assistance
Concept visualization

Use Cases

Creative Design
Concept Art Creation
Quickly generate concept art based on text descriptions
Produces visual works aligned with prompt semantics
Product Design Prototyping
Generate product design sketches for team discussion
Rapid visualization of design concepts
Education & Research
Generative Model Research
Study limitations and biases of diffusion models
Identify and improve socio-cultural biases in models
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