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Stable Diffusion Xl Inpainting 1.0 GGUF

Developed by gpustack
A text-to-image generation model based on Stable Diffusion XL with inpainting capabilities, able to modify image content based on text input and masks.
Downloads 387
Release Time : 12/24/2024

Model Overview

This is a latent text-to-image diffusion model capable of generating realistic images from any text input, with the additional functionality of inpainting images using masks.

Model Features

High-resolution image generation
Supports 1024x1024 high-resolution image generation
Image inpainting capability
Can repair specific areas of an image using masks
Dual text encoder
Uses OpenCLIP-ViT/G and CLIP-ViT/L dual text encoders to enhance text understanding
Quantization support
Provides multiple quantization versions (FP16, Q8_0, Q4_1, Q4_0) to accommodate different hardware requirements

Model Capabilities

Text-to-image generation
Image inpainting
High-resolution image generation
Mask-based content editing

Use Cases

Creative design
Art creation
Generate artworks based on text descriptions
Produces high-quality images matching the description
Image restoration
Repair damaged or unwanted parts of an image
Seamlessly restores specified areas of an image
Educational tools
Visual teaching aids
Quickly generate visual materials for teaching
Generates relevant images based on teaching content
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