flux_lora_v1 Open-Source Model - Freely Achieve Text-to-Image Conversion, Exceptional for Stylized Character Scenes

Flux Lora V1

Developed by stablediffusionapi

LoRA model trained on the FLUX.1-dev base model for text-to-image generation, particularly suitable for stylized character scenes

Image Generation Open Source License:Other #LoRA fine-tuning #FLUX.1 adaptation #Stylized portraits

Downloads 105

Release Time : 2/28/2025

Model Overview

This model is a Low-Rank Adaptation (LoRA) network designed to enhance the text-to-image generation capability of diffusion models, especially excelling in generating stylized character images

Model Features

Stylized character generation

Particularly adept at generating stylized characters in specific environments like bar scenes

Lightweight adaptation

As a LoRA model, it can be flexibly loaded onto different base models without significantly increasing computational burden

Multi-framework compatibility

Supports various mainstream stable diffusion frameworks including ComfyUI and AUTOMATIC1111

Model Capabilities

Text-to-image generation

Stylized character creation

Scene-specific image generation

Use Cases

Creative arts

Bar scene character creation

Generate stylized character images in bar environments

As shown in example images like 'Sexy Sweeney in the bar'

Personalized character design

Create character images with specific personality traits

As shown in example images like 'Lively Roxy'

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Flux Lora V1

Model Overview

Model Features

Model Capabilities

Use Cases

🚀 lora_train_1

🚀 Quick Start

Trigger words

Download model and use it with ComfyUI, AUTOMATIC1111, SD.Next, Invoke AI, etc.

Use it with the 🧨 diffusers library

📄 License