# High-resolution output

Nova Anime XL is a text-to-image generation model based on Stable Diffusion XL technology, specializing in generating photorealistic and anime-style pony images.

A LoRA model based on FLUX.1-dev, specifically designed for generating 3D-style emoji images.

LoRA model based on the Flux.1 development version, specializing in generating high-quality cartoon-style poster images

This is a Diffusers version implementation of the Canny ControlNet checkpoint based on the FLUX.1-dev model, designed for text-to-image generation tasks.

An openpose controlnet model trained for flux-dev, used for controlling image generation through human poses

ControlNet++ is a powerful image generation and editing model that supports multiple control conditions and can generate high-resolution images with visual effects comparable to Midjourney.

A stable diffusion model based on the Kohaku XL Epsilon rev2 version, supporting text-to-image tasks with improved style and character fidelity, and more style options.

Segmind-Vega is a distilled version of Stable Diffusion XL (SDXL), achieving a 70% reduction in model size and a 100% speed improvement while maintaining high-quality text-to-image generation capabilities.

PixArt-LCM is a text-to-image generation model based on the diffusion Transformer, combining the advantages of Pixart-α and LCM. It can quickly generate high-quality images according to text prompts.

A text-to-image generation model based on stable diffusion technology that can generate high-quality images according to text descriptions.

GhostMix is a text-to-image generation model based on stable diffusion technology, focusing on high-quality image generation.

A text-to-image generation model based on stable diffusion technology, capable of generating high-quality images from text descriptions.

A bilingual text-to-image generation model developed based on AltCLIP and Stable Diffusion framework, supporting Chinese and English input