๐ SD3.5-4/8Step-Large-Turbo-HyperRealistic-LoRA
This project offers a LoRA model based on Stable Diffusion 3.5 Large Turbo, aiming to generate hyper - realistic images. It's currently in the training phase and may have some performance issues.
Tags
- text - to - image
- lora
- diffusers
- template:diffusion - lora
- stable - diffusion - 3.5 - large - turbo
- hyper - realistic
- realism
- epic - realism
Widget Examples
- Input: Hyper Realistic, dark haired woman is standing in front of orange light, in glossy transparent plastic futuristic raincoat, short hair, volumetric lighting, cinematic angle, glamour vibrant.
Output: Image
- Input: Photograph of two young men, one with short blonde hair and the other with brown hair, standing in front of mountains and wearing old West - style and tank tops, abs, fit, The image is in a vintage black - and - white photography style, in the style of Richard Avedon. --v 6. 1 --ar 113:128
Output: Image
- Input: A photorealistic portrait of a 22 - year - old Indian woman smiles confidently at the camera, with long, straight black hair and warm brown eyes. She wears a yet casual dress that combines comfort and style, posed in a full - body shot with her hands casually resting at her sides, showcasing a stylish watch on her left wrist. The image features high detail and quality, captured in soft, natural lighting that accentuates her approachable expression as she stands in a modern urban environment.
Output: Image
Base Model
stabilityai/stable - diffusion - 3.5 - large - turbo
Instance Prompt
hyper realistic
License
creativeml - openrail - m
๐ Quick Start
Demo space + Grid Gen Upto 3X2 : https://huggingface.co/spaces/prithivMLmods/SD3.5 - TURBO - REALISM (Not In Use)
โ ๏ธ Important Note
The model is still in the training phase. This is not the final version and may contain artifacts and perform poorly in some cases.
โจ Features
The model can generate hyper - realistic images based on the given prompts, leveraging the power of Stable Diffusion 3.5 Large Turbo and LoRA technology.
๐ฆ Installation
from diffusers import DiffusionPipeline
import torch
device = "cuda" if torch.cuda.is_available() else "cpu"
model_repo_id = "stabilityai/stable-diffusion-3.5-large-turbo"
torch_dtype = torch.bfloat16 if torch.cuda.is_available() else torch.float32
pipe = DiffusionPipeline.from_pretrained(model_repo_id, torch_dtype=torch_dtype)
pipe = pipe.to(device)
pipe.load_lora_weights("prithivMLmods/SD3.5-Large-Turbo-HyperRealistic-LoRA", weight_name="SD3.5-4Step-Large-Turbo-HyperRealistic-LoRA.safetensors")
trigger_word = "hyper realistic"
pipe.fuse_lora(lora_scale=1.0)
๐ Documentation
Model description
prithivMLmods/SD3.5 - Large - Turbo - HyperRealistic - LoRA
Image Processing Parameters
Parameter |
Value |
Parameter |
Value |
LR Scheduler |
constant |
Noise Offset |
0.03 |
Optimizer |
AdamW |
Multires Noise Discount |
0.1 |
Network Dim |
64 |
Multires Noise Iterations |
10 |
Network Alpha |
32 |
Repeat & Steps |
20 & 2.2k |
Epoch |
15 |
Save Every N Epochs |
1 |
Labeling: florence2 - en(natural language & English)
Total Images Used for Training : 30
Sample Image


App File Structure
/project - root/
โโโ .gitattributes
โโโ README.md
โโโ app.py
โโโ pythonproject.py
Trigger words ๐งจ
โ ๏ธ Important Note
You should use hyper realistic
to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
[Download](/prithivMLmods/SD3.5 - 4Step - Large - Turbo - HyperRealistic - LoRA/tree/main) them in the Files & versions tab.