# Text-to-Image

The GGUF quantized version of the Chroma model, supporting text-to-image conversion and image generation.

This is a DreamBooth LoRA model based on HiDream-ai/HiDream-I1-Full, specializing in generating yarn art-style dog images.

HiDream-I1-Fast is a GGUF-quantized text-to-image generation model designed for fast image generation tasks.

Hierarchical Distillation (SwD) is a novel framework that accelerates diffusion models (DMs) by progressively increasing spatial resolution during the generation process.

This is a LoRA model trained based on FLUX.1-dev, specifically designed for generating facial images of a particular style.

doa is a text-to-image generation project created based on the SDXL model, developed by reijlita, which can achieve high-quality text-to-image conversion.

Flux Handwriting is a LoRA adapter based on the FLUX.1-dev model, specializing in generating various styles of handwritten text images.

A text-to-image model based on Stable Diffusion XL, supporting ControlNet technology, suitable for high-quality image generation tasks.

Flux LoRA model for creating fantasy illustrations, especially suitable for generating images in WHMSCPE001 style

GGUF quantized version of FLUX.1-dev, focused on text-to-image and image generation tasks

GGUF quantized version of the FLUX.1-dev model, focused on text-to-image and image generation tasks

Lumina-mGPT is a family of multimodal autoregressive models excelling in various vision and language tasks, particularly in generating flexible and realistic images from text descriptions.

Lumina-mGPT is a family of multimodal autoregressive models, excelling in generating flexible and realistic images from text descriptions and capable of performing various vision and language tasks.

Lumina-mGPT is a family of multimodal autoregressive models, excelling in generating flexible and realistic images from text descriptions, and capable of performing various vision and language tasks.

Lumina-mGPT is a series of multimodal autoregressive models, excelling in generating flexible and realistic images from text descriptions.

This is a text-to-image model based on stable diffusion technology, capable of generating high-quality images from input text descriptions.

This is a LoRA adapter weight for Stable Diffusion XL (SDXL), specifically designed for generating advertising images with a particular style.

Stable-diffusion-GGUF is a text-to-image generation model that offers multiple quantization versions and is suitable for fields such as art creation.

This is a LoRA model based on Stable Diffusion XL, specifically designed for generating glamour shot-style images.

This is a LoRA-adapted weight model based on Stable Diffusion XL (SDXL), specifically designed to generate images with a distinctive seaweed painting style.

This model can generate high-quality images based on input text descriptions, suitable for various scenarios such as creative design and content creation.

The Latent Consistency Model is a text-to-image generation model distilled from a fine-tuned version of Stable-Diffusion v1-5's Dreamshaper v7, capable of producing high-quality images in extremely short inference times.

The Stable Diffusion Model is a text-to-image diffusion model capable of generating high-quality images based on textual descriptions.

A text-to-image generation model based on diffusion models, capable of producing high-quality images from textual descriptions.

This is a text-to-image generation model fine-tuned based on Stable Diffusion v1-5, specifically designed for generating emoji-style images.

A text-to-image model fine-tuned from Stable Diffusion, specialized in generating visually stunning robot images

A Stable Diffusion model fine-tuned on a Russian emoji dataset, specialized in generating emoji images

A Stable Diffusion model trained on the art style of 'Bloodborne', capable of generating images with the game's unique artistic style

This is a Stable Diffusion model fine-tuned on El Risitas images, designed to generate images in the style of El Risitas.

Stable Diffusion model trained using the Diffusers library, optimized with a One Piece anime-related dataset