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Stable Diffusion V1 5 GGUF

Developed by gpustack
Stable Diffusion v1-5 is a latent text-to-image diffusion model capable of generating realistic images from text inputs.
Downloads 996
Release Time : 11/8/2024

Model Overview

Stable Diffusion is a latent diffusion model that can generate and modify images based on text prompts. It uses a fixed pre-trained text encoder (CLIP ViT-L/14) and is suitable for research purposes and creative applications.

Model Features

High-Quality Image Generation
Capable of generating high-quality, realistic images from text prompts.
Classifier-Free Guidance Sampling
Improves sampling by partially discarding text conditions.
Multi-Framework Support
Supports various frameworks such as Diffusers library, ComfyUI, and Automatic1111.

Model Capabilities

Text-to-Image Generation
Image Modification
Artistic Creation

Use Cases

Art & Design
Generate Artwork
Create unique artworks based on creative text prompts.
Generates high-quality artistic images
Education & Research
Generative Model Research
Used to study the limitations and biases of generative models.
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